Stable Diffusion XL-base-1.0

Stable Diffusion XL-base-1.0 is the base* text-to-image generation model, improved upon previous versions of Stable Diffusion models. It is designed to generate images at 1024x1024 pixels. It is also not recommended to choose an image size smaller than 512x512 pixels.

*SDXL consists of an ensemble of experts pipeline for latent diffusion: In a first step, the base model is used to generate (noisy) latents, which are then further processed with a refinement model specialized for the final denoising steps. Note that the base model can be used as a standalone module.

stabilityai/stable-diffusion-xl-base-1.0 // https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0

Announce Date: 26.07.2023
Parameters: 3.5B
Context: 77
VRAM requirements: 1.6 GB using 4 bits quantization
Developer: Stability AI
License: OpenRail++

Use our pre-built public endpoints to test inference and explore Stable Diffusion XL-base-1.0 capabilities.

Rent your own physically dedicated instance with hourly or long-term monthly billing.

We recommend deploying private instances in the following scenarios:

• maximize endpoint performance,
• enable full context for long sequences,
• ensure top-tier security for data processing in an isolated, dedicated environment,
• use custom weights, such as fine-tuned models or LoRA adapters.

Contact our dedicated neural networks support team at nn@immers.cloud or send your request to the sales department at sale@immers.cloud.