The Stable Diffusion v2-1 is a diffusion-based text-to-image generation model. It is fine-tuned from the Stable Diffusion v2 checkpoint. Generates high-resolution images (up to 768x768) and supports text-guided generation.
Stable Diffusion x4 upscaler - is a diffusion model for 4x image resolution upscaling based on text prompts. It takes as input a low-resolution image and a text prompt, along with the noise_level parameter.
Stable Diffusion v1.5 is a diffusion model for image generation based on textual prompts. The model was initialized with the weights of the previous version, Stable Diffusion v1.2, and subsequently fine-tuned. It supports image generation at 512×512 pixels and image modification via textual prompts.