Stable Diffusion XL-base-1.0 is the base text-to-image generation model, improved upon previous versions of Stable Diffusion models. It is designed to generate images at 1024x1024 pixels. It is also not recommended to choose an image size smaller than 512x512 pixels.
Refiner Model specializes in the final stages of noise reduction and enhances the visual accuracy of images generated by the base model.
Kandinsky 2.2 is a free Russian model for image generation developed by Sber AI. It operates based on a diffusion model: initially adding noise to images it was trained on, then restoring them through a reverse diffusion process to create new unique images.
The Stable Diffusion v2-1 is a diffusion-based text-to-image generation model. It is fine-tuned from the Stable Diffusion v2 checkpoint. Generates high-resolution images (up to 768x768) and supports text-guided generation.
Stable Diffusion x4 upscaler - is a diffusion model for 4x image resolution upscaling based on text prompts. It takes as input a low-resolution image and a text prompt, along with the noise_level parameter.
Stable Diffusion v1.5 is a diffusion model for image generation based on textual prompts. The model was initialized with the weights of the previous version, Stable Diffusion v1.2, and subsequently fine-tuned. It supports image generation at 512×512 pixels and image modification via textual prompts.